mnist diffusion
时间: 2023-08-04 21:09:10 浏览: 187
MNIST是一个经典的手写数字识别数据集,其中包含了来自0到9的手写数字图像。Diffusion是一种基于梯度的图像处理技术,它可以用于增强图像的细节和边缘,使图像更加清晰和有结构。在MNIST数据集中应用Diffusion技术可能会产生一些有趣的结果,例如增强数字的边缘和细节,使它们更易于识别和分类。这种技术可以帮助改进基于MNIST数据集的手写数字识别算法。
相关问题
diffusion model 生成mnist
### 使用 Diffusion Model 生成 MNIST 数据集图像
为了使用扩散模型 (Diffusion Model, DM) 来生成 MNIST 数据集中的手写数字图像,可以遵循以下方法。此过程涉及两个主要阶段:前向传播(即逐渐增加噪声)以及反向传播(即去噪)。具体来说:
#### 前向传播(加噪)
在这一过程中,原始图像被逐步加入高斯白噪音直到变得完全随机化。这一步骤定义了一系列由 $\beta_t$ 控制的方差参数下的马尔可夫链状态转移[^1]。
#### 后向传播(去噪)
该部分旨在逆转上述加噪流程,从纯噪声恢复到清晰图像。网络学习预测每一步中应减去多少噪声以接近初始输入。整个训练目标是最小化均方误差损失函数,在这里用来衡量真实分布 $q(\mathbf{x}_t|\mathbf{x}_{t-1})$ 和估计分布 $p_\theta (\mathbf{x}_{t-1} | \mathbf{x}_t)$之间的差异。
下面是基于PyTorch的一个简化版实现例子,用于演示如何构建并应用一个简单的扩散模型来进行MNIST数据集上的图像生成任务:
```python
import torch
from torch import nn
import torch.nn.functional as F
from torchvision.datasets import MNIST
from torchvision.transforms import Compose, ToTensor, Lambda
from torch.utils.data import DataLoader
class SimpleUnet(nn.Module):
def __init__(self):
super().__init__()
self.encoder = nn.Sequential(
nn.Conv2d(1, 64, kernel_size=3),
nn.ReLU(),
nn.BatchNorm2d(64),
nn.Conv2d(64, 128, kernel_size=3),
nn.MaxPool2d(kernel_size=2),
nn.ReLU(),
nn.BatchNorm2d(128),
nn.Conv2d(128, 256, kernel_size=3),
nn.MaxPool2d(kernel_size=2),
nn.ReLU()
)
self.decoder = nn.Sequential(
nn.Upsample(scale_factor=2),
nn.ConvTranspose2d(256, 128, kernel_size=3),
nn.ReLU(),
nn.BatchNorm2d(128),
nn.Upsample(scale_factor=2),
nn.ConvTranspose2d(128, 64, kernel_size=3),
nn.ReLU(),
nn.BatchNorm2d(64),
nn.ConvTranspose2d(64, 1, kernel_size=3),
nn.Sigmoid()
)
def forward(self, x):
encoded_x = self.encoder(x)
decoded_x = self.decoder(encoded_x)
return decoded_x
def linear_beta_schedule(timesteps, start=0.0001, end=0.02):
return torch.linspace(start, end, timesteps)
if __name__ == "__main__":
device = "cuda" if torch.cuda.is_available() else "cpu"
T = 300
betas = linear_beta_schedule(T).to(device)
alphas = 1 - betas
alpha_bars = torch.cumprod(alphas, axis=0)
unet_model = SimpleUnet().to(device)
optimizer = torch.optim.Adam(unet_model.parameters(), lr=1e-3)
transform = Compose([ToTensor(), Lambda(lambda x: (x * 2) - 1)])
dataset = MNIST(root='./data', train=True, download=True, transform=transform)
dataloader = DataLoader(dataset, batch_size=128, shuffle=True)
epochs = 10
for epoch in range(epochs):
for step, (images, labels) in enumerate(dataloader):
images = images.to(device)
# Sample a random timestep from the distribution.
ts = torch.randint(low=0, high=T, size=(len(images),)).long().to(device)
noise = torch.randn_like(images).to(device)
noisy_images = (
sqrt_alpha_bar[ts][:, None, None, None].sqrt()*images +
(1-alpha_bars[ts])[:,None,None,None].sqrt()*noise
).clamp(-1., 1.)
predicted_noise = unet_model(noisy_images, ts)
loss = F.mse_loss(predicted_noise, noise)
optimizer.zero_grad()
loss.backward()
optimizer.step()
if not step % 100:
print(f'Epoch {epoch}, Step {step}: Loss={loss.item()}')
```
这段代码展示了怎样创建一个简易版本的U-net架构作为基础模型,并利用线性的$\beta$调度策略来控制噪声水平变化。此外还包含了完整的训练循环逻辑,其中涉及到采样时间步$t\in[T]$、计算带噪图片$x^{(t)}=\sqrt{\bar{\alpha}}_tx+(1-\bar{\alpha}_t)\epsilon,\;\epsilon∼N(0,I)$以及最小化预测噪声与实际添加噪声间的MSE距离[^2].
生成一个基于pytorch深度学习框架,数据集采用fashion-mnist,模型使用Diffusion模型,添加数据可视化的,用Python语言写的完整代码
以下是一个基于PyTorch深度学习框架,使用Diffusion模型处理Fashion-MNIST数据集,并添加数据可视化的完整Python代码:
```python
import torch
import torch.nn as nn
import torch.optim as optim
from torchvision import datasets, transforms
import matplotlib.pyplot as plt
import numpy as np
# 设置设备
device = torch.device("cuda" if torch.cuda.is_available() else "cpu")
# 超参数
batch_size = 128
epochs = 20
learning_rate = 1e-3
image_size = 28
channels = 1
n_steps = 1000
beta_start = 1e-4
beta_end = 0.02
# 数据预处理
transform = transforms.Compose([
transforms.ToTensor(),
transforms.Normalize((0.5,), (0.5,))
])
# 加载数据集
train_dataset = datasets.FashionMNIST(root='./data', train=True, transform=transform, download=True)
train_loader = torch.utils.data.DataLoader(train_dataset, batch_size=batch_size, shuffle=True)
# 定义扩散模型
class DiffusionModel(nn.Module):
def __init__(self):
super(DiffusionModel, self).__init__()
self.model = nn.Sequential(
nn.Conv2d(1, 64, kernel_size=3, stride=1, padding=1),
nn.ReLU(),
nn.Conv2d(64, 64, kernel_size=3, stride=2, padding=1),
nn.ReLU(),
nn.Conv2d(64, 128, kernel_size=3, stride=2, padding=1),
nn.ReLU(),
nn.Flatten(),
nn.Linear(128 * 7 * 7, 1024),
nn.ReLU(),
nn.Linear(1024, 128 * 7 * 7),
nn.ReLU(),
nn.Unflatten(1, (128, 7, 7)),
nn.ConvTranspose2d(128, 64, kernel_size=3, stride=2, padding=1, output_padding=1),
nn.ReLU(),
nn.ConvTranspose2d(64, 64, kernel_size=3, stride=2, padding=1, output_padding=1),
nn.ReLU(),
nn.Conv2d(64, 1, kernel_size=3, stride=1, padding=1),
nn.Tanh()
)
def forward(self, x):
return self.model(x)
# 初始化模型、损失函数和优化器
model = DiffusionModel().to(device)
criterion = nn.MSELoss()
optimizer = optim.Adam(model.parameters(), lr=learning_rate)
# 训练循环
for epoch in range(epochs):
for i, (images, _) in enumerate(train_loader):
images = images.to(device)
# 生成时间步
t = torch.randint(0, n_steps, (batch_size,), device=device).long()
# 计算beta
betas = torch.linspace(beta_start, beta_end, n_steps).to(device)
alpha = 1 - betas
alpha_hat = torch.cumprod(alpha, dim=0)
# 获取噪声
noise = torch.randn_like(images).to(device)
# 计算noisy images
sqrt_alpha_hat = torch.sqrt(alpha_hat[t].reshape(-1, 1, 1, 1))
sqrt_one_minus_alpha_hat = torch.sqrt(1 - alpha_hat[t].reshape(-1, 1, 1, 1))
noisy_images = sqrt_alpha_hat * images + sqrt_one_minus_alpha_hat * noise
# 前向传播
outputs = model(noisy_images)
loss = criterion(outputs, noise)
# 反向传播和优化
optimizer.zero_grad()
loss.backward()
optimizer.step()
print(f"Epoch [{epoch+1}/{epochs}], Loss: {loss.item():.4f}")
# 可视化结果
with torch.no_grad():
# 生成随机噪声
sample = torch.randn(64, channels, image_size, image_size).to(device)
# 反向扩散过程
for i in range(n_steps):
betas_t = beta_end * torch.ones(batch_size, device=device)
alpha_t = 1 - betas_t
alpha_hat_t = torch.cumprod(alpha_t, dim=0)[-1]
sqrt_recip_alpha_t = torch.sqrt(1 / alpha_t)
sqrt_one_minus_alpha_hat_t = torch.sqrt(1 - alpha_hat_t)
noise = torch.randn_like(sample).to(device)
# 模型预测
model_input = sample
model_output = model(model_input)
# 更新样本
sample = sqrt_recip_alpha_t * (model_input - (1 - alpha_t) / torch.sqrt(1 - alpha_hat_t) * model_output) + torch.sqrt(betas_t) * noise
# 显示生成的图片
grid = utils.make_grid(sample.cpu(), nrow=8)
plt.figure(figsize=(8,8))
plt.imshow(grid.permute(1, 2, 0).squeeze(), cmap='gray')
plt.axis('off')
plt.show()
```
这段代码实现了一个基本的Diffusion模型,用于处理Fashion-MNIST数据集,并包含了数据可视化功能。
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